A study of language models, especially those designed to create images and video, might suggest that copyright infringements by diffusion-based models might be impossible to prove.
The authors say “locating a part of the training data that can be held responsible for a generated sample, can become impossible if a model is trained on a sufficiently large corpus of data.”
In other words, a large training set is all that is needed to render copyright traceable to any single source impossible.
The study methodology relies on omitting an attributed unit from the training set (one image or all images by the same creator), the induced counterfactual sample cannot be directly attributed to that same unit.

source: Nature Communications
In the above example, a generated image created from a smaller data set is more attributable to a single training image.
A large data set, on the other hand, creates a new image that shows no trace of a single training image.
Model / family | Main generation | Diffusion? | Basic architecture / approach |
GPT-4 / GPT-4o / GPT-5 family | Text, multimodal | No* | Primarily autoregressive Transformer for language |
Claude 3/4 family | Text, code, multimodal | No* | Autoregressive Transformer |
Gemini 1–4 family | Text, multimodal | No* | Primarily autoregressive Transformer |
Llama family | Text, code | No | Autoregressive Transformer |
Mistral / Mixtral family | Text, code | No | Autoregressive Transformer, including MoE |
DeepSeek family | Text, code, reasoning | No | Autoregressive Transformer/MoE |
Gemini Diffusion | Text/code | Yes | Text diffusion Transformer; generates blocks of tokens and iteratively refines them (Google DeepMind) |
DiffusionGemma | Text/code | Yes | Non-sequential diffusion Transformer based on Gemma research (Google DeepMind) |
Stable Diffusion 1.x/2.x/XL | Images | Yes | Latent diffusion, historically U-Net-based |
Stable Diffusion 3/3.5 | Images | Yes | Diffusion Transformer / MM-DiT |
FLUX.1 / FLUX family | Images | Yes | Diffusion/flow-based Transformer |
DALL·E 2 | Images | Yes | Diffusion |
DALL·E 3 | Images | Yes | Diffusion-based image generation |
Imagen 2/3/4 | Images | Yes | Diffusion-based image generation; Google explicitly describes Imagen as a diffusion model (Google DeepMind) |
Midjourney | Images | Yes* | Widely understood to use diffusion/related denoising techniques, though Midjourney does not disclose its architecture in detail |
Adobe Firefly | Images/video | Yes* | Diffusion-based generation among its models; architecture varies by model |
Sora | Video | Yes | Diffusion Transformer (DiT); starts from noisy video patches and denoises them (OpenAI) |
Sora 2 | Video | Yes* | Diffusion-based video generation; details of current implementation are less fully disclosed |
Veo 2/3/3.1 | Video | Yes* | Diffusion/Transformer-based video generation; Google describes Veo as drawing on its diffusion-model research (Google DeepMind) |
Runway Gen-2/3 | Video | Yes* | Diffusion-based video generation |
Kling | Video | Yes* | Diffusion/Transformer-based video generation |
Luma Dream Machine | Video | Yes* | Diffusion/Transformer-based video generation |
MusicGen | Music/audio | No | Autoregressive Transformer over discrete audio tokens |
AudioCraft / related models | Audio/music | Mixed | Different models use different architectures |
Genie 2 | Interactive worlds/video | Yes | Autoregressive latent diffusion—an interesting hybrid (Google DeepMind) |
“Our findings provide a compelling case that attribution, the task of locating a unit of data within the training set that can be held responsible for a generated sample, is practically impossible on contemporary generative diffusion models,” say authors Zheng Dai and David Gifford.
No comments:
Post a Comment